Created
September 16, 2022 01:38
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Creating images with Stable Diffusion to find a good seed to go with prompt
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import torch | |
from diffusers import StableDiffusionPipeline | |
from torch import autocast | |
import random | |
import matplotlib.pyplot as plt | |
import os | |
prompts = [ | |
"1965 Porsche 911", | |
"1975 Porsche 911", | |
"1985 Porsche 911", | |
"1995 Porsche 911", | |
"2005 Porsche 911 front", | |
"2015 Porsche 911", | |
"2020 Porsche 911", | |
"2020 Porsche 911 GT3 RS"] | |
# make sure you're logged in with `huggingface-cli login` | |
pipe = StableDiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", | |
revision="fp16", | |
torch_dtype=torch.float16, | |
use_auth_token=True) | |
pipe = pipe.to("cuda") | |
def infer(prompt, num_inference_steps=50, | |
samples=5, seed=1024, guidance_scale=7.5, | |
width=512, height=512): | |
generator = torch.Generator("cuda").manual_seed(seed) | |
w = width//8*8 | |
h = height//8*8 | |
with autocast("cuda"): | |
image = pipe(prompt, guidance_scale=7.5, | |
generator=generator, width=w, height=h, | |
num_inference_steps=num_inference_steps)["sample"][0] | |
return image | |
for p in prompts: | |
prompt_orig = p | |
if not os.path.exists("imagery"): | |
os.mkdir("imagery") | |
if not os.path.exists(f"imagery/{prompt_orig}"): | |
os.mkdir(f"imagery/{prompt_orig}/") | |
HM = 200 | |
for i in range(HM): | |
print(f"{i+1}/{HM}") | |
prompt_to_use = prompt_orig | |
seed = random.randint(0, 10000) | |
print(seed) | |
image = infer(prompt_to_use, num_inference_steps=75, | |
seed=seed, width=512, height=512) | |
image.save(f"imagery/{prompt_to_use}/{seed}.png") |
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Hi was looking for something similar so thanks for this. Does this work with things like text embeddings and LORA or would I need to adjust the code?